# Artistic Creation

A diffusion model trained with LoRA on storyboard sketch data, designed to generate images in a storyboard sketch style.

SOAP is a model focused on generating stylized animated portraits, capable of producing 3D animation effects from input images.

LoRA adapter trained on FLUX.1-dev model for text-to-image generation

Stable Diffusion is a text-to-image diffusion model capable of generating high-quality images from text descriptions.

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from text inputs.

Stable Diffusion v1-5 is a latent text-to-image diffusion model capable of generating realistic images from text inputs.

UniPortrait is a text-to-image generation model based on the diffusers library, capable of generating high-quality portrait images from text descriptions.

Image-to-image conversion model based on black-forest-labs/FLUX.1-dev

A Stable Diffusion XL model specialized in illustration generation, developed based on Kohaku XL Beta 5 and trained by the OnomaAI research team

FLUX.1-dev is a text-to-image generation model based on the diffusers library, focusing on high-quality image synthesis.

Flux LoRA model for creating fantasy illustrations, especially suitable for generating images in WHMSCPE001 style

A latent diffusion model for text-to-image generation, supporting 512x512 resolution image generation

A LoRA model based on Stable Diffusion, designed to generate fantastical creatures resembling a hybrid of kangaroos and deer, with flowering antlers and tails that glow in the dark.

A latent diffusion model for text-to-image generation that can create realistic images from arbitrary text inputs

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from any text input.

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from text inputs.

A text-to-image generation model based on latent diffusion model with image inpainting capability

Stable Diffusion 1.5 is a high-resolution image generation model based on latent diffusion models, capable of generating high-quality images from text descriptions.

A high-resolution image generation model based on latent diffusion models, capable of producing high-quality artistic images from text descriptions.

A high-resolution image generation model based on latent diffusion models, supporting text-to-image and image-to-image translation.

GGUF quantized version of Stable Diffusion 2.1, supporting multiple quantization precisions, suitable for text-to-image generation tasks.

A text-to-image generation model based on StableDiffusionXLPipeline, focused on artistic creation.

Midjourney is the world's most realistic and powerful AI image generator, marking the first release of its V7 version.

A text-to-image generation model trained by rd690 based on NxtWave's 'Build Your Own Gen AI Model' course, specializing in animal-themed image generation.

This is a LoRA-adapted weight model based on Stable Diffusion XL (SDXL), specifically designed to generate images with a distinctive seaweed painting style.

This is a LoRA model based on Stable Diffusion XL, specializing in generating images with an 'entrancing' style.

This is an unconditional image generation model based on PyTorch and Diffusers library, specifically designed for generating anime-style face images.

A distilled consistency adapter designed for segmind/SSD-1B, reducing inference steps to just 2-8 steps

A latent consistency model based on Stable Diffusion XL, reducing inference steps to 2-8

The Stable Diffusion Model is a text-to-image diffusion model capable of generating high-quality images based on textual descriptions.

A LoRA model based on Stable Diffusion XL, specifically designed for generating pencil sketch-style images.

A text-to-image generation model based on diffusion models, capable of producing high-quality images from textual descriptions.

SD-XL 0.9-refiner is a latent diffusion model developed by Stability AI, designed for high-quality image optimization and requires pairing with a base model

ControlNet v1.1 is a neural network architecture based on Stable Diffusion, designed to control diffusion models through image inpainting conditions.

A text-to-image generation model based on latent diffusion architecture with enhanced image inpainting capabilities using masks

A fine-tuned version of Stable Diffusion 2.1 that supports generating and modifying images through text prompts or noisy CLIP image embeddings

A diffusion-based text-to-image generation model that supports generating and editing high-quality images through text prompts

A fine-tuned image generation model based on Stable Diffusion 1.5, specializing in high-quality photorealistic style images.

A latent diffusion model for generating high-quality images from text prompts, supporting 512x512 resolution image generation

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from any text input.